xiefan-guo/initno

[CVPR 2024] InitNO: Boosting Text-to-Image Diffusion Models via Initial Noise Optimization

29
/ 100
Experimental

This project helps anyone generating images from text descriptions using Stable Diffusion models to get higher-quality, more accurate results. It takes your existing text prompts and the initial noise used by the model, then optimizes that noise to ensure the generated image closely matches your textual intent. This is for users who already work with text-to-image diffusion models and want to improve the reliability and fidelity of their outputs.

No commits in the last 6 months.

Use this if you are a digital artist, content creator, or researcher frequently generating images with Stable Diffusion and often find the initial outputs don't quite match your text prompts due to 'bad' starting noise.

Not ideal if you are looking for a new text-to-image model from scratch or are not familiar with existing diffusion model workflows.

AI-art-generation digital-content-creation image-synthesis text-to-image generative-AI-fine-tuning
Stale 6m No Package No Dependents
Maintenance 0 / 25
Adoption 9 / 25
Maturity 16 / 25
Community 4 / 25

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77

Forks

2

Language

Python

License

Apache-2.0

Last pushed

Jun 07, 2024

Commits (30d)

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